State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.
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SKILL.md
Quellanweisungen · Schreibgeschützte Vorschau
name
stable-diffusion-image-generation
description
State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.
Text Prompt → Text Encoder → Text Embeddings
↓
Random Noise → [Denoising Loop] ← Scheduler
↓
Predicted Noise
↓
VAE Decoder → Final Image
Core concepts
Pipelines
Pipelines orchestrate complete workflows:
Pipeline
Purpose
StableDiffusionPipeline
Text-to-image (SD 1.x/2.x)
StableDiffusionXLPipeline
Text-to-image (SDXL)
StableDiffusion3Pipeline
Text-to-image (SD 3.0)
FluxPipeline
Text-to-image (Flux models)
StableDiffusionImg2ImgPipeline
Image-to-image
StableDiffusionInpaintPipeline
Inpainting
Schedulers
Schedulers control the denoising process:
Scheduler
Steps
Quality
Use Case
EulerDiscreteScheduler
20-50
Good
Default choice
EulerAncestralDiscreteScheduler
20-50
Good
More variation
DPMSolverMultistepScheduler
15-25
Excellent
Fast, high quality
DDIMScheduler
50-100
Good
Deterministic
LCMScheduler
4-8
Good
Very fast
UniPCMultistepScheduler
15-25
Excellent
Fast convergence
Swapping schedulers
from diffusers import DPMSolverMultistepScheduler
# Swap for faster generation
pipe.scheduler = DPMSolverMultistepScheduler.from_config(
pipe.scheduler.config
)
# Now generate with fewer steps
image = pipe(prompt, num_inference_steps=20).images[0]
Generation parameters
Key parameters
Parameter
Default
Description
prompt
Required
Text description of desired image
negative_prompt
None
What to avoid in the image
num_inference_steps
50
Denoising steps (more = better quality)
guidance_scale
7.5
Prompt adherence (7-12 typical)
height, width
512/1024
Output dimensions (multiples of 8)
generator
None
Torch generator for reproducibility
num_images_per_prompt
1
Batch size
Reproducible generation
import torch
generator = torch.Generator(device="cuda").manual_seed(42)
image = pipe(
prompt="A cat wearing a top hat",
generator=generator,
num_inference_steps=50
).images[0]
Negative prompts
image = pipe(
prompt="Professional photo of a dog in a garden",
negative_prompt="blurry, low quality, distorted, ugly, bad anatomy",
guidance_scale=7.5
).images[0]
Image-to-image
Transform existing images with text guidance:
from diffusers import AutoPipelineForImage2Image
from PIL import Image
pipe = AutoPipelineForImage2Image.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16
).to("cuda")
init_image = Image.open("input.jpg").resize((512, 512))
image = pipe(
prompt="A watercolor painting of the scene",
image=init_image,
strength=0.75, # How much to transform (0-1)
num_inference_steps=50
).images[0]
Inpainting
Fill masked regions:
from diffusers import AutoPipelineForInpainting
from PIL import Image
pipe = AutoPipelineForInpainting.from_pretrained(
"runwayml/stable-diffusion-inpainting",
torch_dtype=torch.float16
).to("cuda")
image = Image.open("photo.jpg")
mask = Image.open("mask.png") # White = inpaint region
result = pipe(
prompt="A red car parked on the street",
image=image,
mask_image=mask,
num_inference_steps=50
).images[0]
ControlNet
Add spatial conditioning for precise control:
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
import torch
# Load ControlNet for edge conditioning
controlnet = ControlNetModel.from_pretrained(
"lllyasviel/control_v11p_sd15_canny",
torch_dtype=torch.float16
)
pipe = StableDiffusionControlNetPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
controlnet=controlnet,
torch_dtype=torch.float16
).to("cuda")
# Use Canny edge image as control
control_image = get_canny_image(input_image)
image = pipe(
prompt="A beautiful house in the style of Van Gogh",
image=control_image,
num_inference_steps=30
).images[0]
Available ControlNets
ControlNet
Input Type
Use Case
canny
Edge maps
Preserve structure
openpose
Pose skeletons
Human poses
depth
Depth maps
3D-aware generation
normal
Normal maps
Surface details
mlsd
Line segments
Architectural lines
scribble
Rough sketches
Sketch-to-image
LoRA adapters
Load fine-tuned style adapters:
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16
).to("cuda")
# Load LoRA weights
pipe.load_lora_weights("path/to/lora", weight_name="style.safetensors")
# Generate with LoRA style
image = pipe("A portrait in the trained style").images[0]
# Adjust LoRA strength
pipe.fuse_lora(lora_scale=0.8)
# Unload LoRA
pipe.unload_lora_weights()
Multiple LoRAs
# Load multiple LoRAs
pipe.load_lora_weights("lora1", adapter_name="style")
pipe.load_lora_weights("lora2", adapter_name="character")
# Set weights for each
pipe.set_adapters(["style", "character"], adapter_weights=[0.7, 0.5])
image = pipe("A portrait").images[0]
Memory optimization
Enable CPU offloading
# Model CPU offload - moves models to CPU when not in use
pipe.enable_model_cpu_offload()
# Sequential CPU offload - more aggressive, slower
pipe.enable_sequential_cpu_offload()
Attention slicing
# Reduce memory by computing attention in chunks
pipe.enable_attention_slicing()
# Or specific chunk size
pipe.enable_attention_slicing("max")
from diffusers import UNet2DConditionModel, AutoencoderKL
# Load custom VAE
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse")
# Use with pipeline
pipe = DiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
vae=vae,
torch_dtype=torch.float16
)
Batch generation
Generate multiple images efficiently:
# Multiple prompts
prompts = [
"A cat playing piano",
"A dog reading a book",
"A bird painting a picture"
]
images = pipe(prompts, num_inference_steps=30).images
# Multiple images per prompt
images = pipe(
"A beautiful sunset",
num_images_per_prompt=4,
num_inference_steps=30
).images
Common workflows
Workflow 1: High-quality generation
from diffusers import StableDiffusionXLPipeline, DPMSolverMultistepScheduler
import torch
# 1. Load SDXL with optimizations
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16"
)
pipe.to("cuda")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe.enable_model_cpu_offload()
# 2. Generate with quality settings
image = pipe(
prompt="A majestic lion in the savanna, golden hour lighting, 8k, detailed fur",
negative_prompt="blurry, low quality, cartoon, anime, sketch",
num_inference_steps=30,
guidance_scale=7.5,
height=1024,
width=1024
).images[0]
Workflow 2: Fast prototyping
from diffusers import AutoPipelineForText2Image, LCMScheduler
import torch
# Use LCM for 4-8 step generation
pipe = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16
).to("cuda")
# Load LCM LoRA for fast generation
pipe.load_lora_weights("latent-consistency/lcm-lora-sdxl")
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
pipe.fuse_lora()
# Generate in ~1 second
image = pipe(
"A beautiful landscape",
num_inference_steps=4,
guidance_scale=1.0
).images[0]
Common issues
CUDA out of memory:
# Enable memory optimizations
pipe.enable_model_cpu_offload()
pipe.enable_attention_slicing()
pipe.enable_vae_slicing()
# Or use lower precision
pipe = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
Black/noise images:
# Check VAE configuration# Use safety checker bypass if needed
pipe.safety_checker = None# Ensure proper dtype consistency
pipe = pipe.to(dtype=torch.float16)
