stable-diffusion-image-generation

# Stable Diffusion Image Generation Comprehensive guide to generating images with Stable Diffusion using the HuggingFace Diffusers library. ## When to use Stable Diffusion **Use Stable Diffusion when:** - Generating images from text descriptions - Performing image-to-image translation (style transfer, enhancement) - Inpainting (filling in masked regions) - Outpainting (extending images beyond boundaries) - Creating variations of existing images - Building custom image generation workflows **Key features:** - **Text-to-Image**: Generate images from natural language prompts - **Image-to-Image**: Transform existing images with text guidance - **Inpainting**: Fill masked regions with context-aware content - **ControlNet**: Add spatial conditioning (edges, poses, depth) - **LoRA Support**: Efficient fine-tuning and style adaptation - **Multiple Models**: SD 1.5, SDXL, SD 3.0, Flux support **Use alternatives instead:** - **DALL-E 3**: For API-based generation without GPU - **Midjourney**: For artistic, stylized outputs - **Imagen**: For Google Cloud integration - **Leonardo.ai**: For web-based creative workflows ## Quick start ### Installation ```bash pip install diffusers transformers accelerate torch pip install xformers # Optional: memory-efficient attention ``` ### Basic text-to-image ```python from diffusers import DiffusionPipeline import torch # Load pipeline (auto-detects model type) pipe = DiffusionPipeline.from_pretrained( "stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16 ) pipe.to("cuda") # Generate image image = pipe( "A serene mountain landscape at sunset, highly detailed", num_inference_steps=50, guidance_scale=7.5 ).images[0] image.save("output.png") ``` ### Using SDXL (higher quality) ```python from diffusers import AutoPipelineForText2Image import torch pipe = AutoPipelineForText2Image.from_pretrained( "stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16" ) pipe.to("cuda") # Enable memory optimization pipe.enable_model_cpu_offload() image = pipe( prompt="A futuristic city with flying cars, cinematic lighting", height=1024, width=1024, num_inference_steps=30 ).images[0] ``` ## Architecture overview ### Three-pillar design Diffusers is built around three core components: ``` Pipeline (orchestration) ├── Model (neural networks) │ ├── UNet / Transformer (noise prediction) │ ├── VAE (latent encoding/decoding) │ └── Text Encoder (CLIP/T5) └── Scheduler (denoising algorithm) ``` ### Pipeline inference flow ``` Text Prompt → Text Encoder → Text Embeddings ↓ Random Noise → [Denoising Loop] ← Scheduler ↓ Predicted Noise ↓ VAE Decoder → Final Image ``` ## Core concepts ### Pipelines Pipelines orchestrate complete workflows: | Pipeline | Purpose | |----------|---------| | `StableDiffusionPipeline` | Text-to-image (SD 1.x/2.x) | | `StableDiffusionXLPipeline` | Text-to-image (SDXL) | | `StableDiffusion3Pipeline` | Text-to-image (SD 3.0) | | `FluxPipeline` | Text-to-image (Flux models) | | `StableDiffusionImg2ImgPipeline` | Image-to-image | | `StableDiffusionInpaintPipeline` | Inpainting | ### Schedulers Schedulers control the denoising process: | Scheduler | Steps | Quality | Use Case | |-----------|-------|---------|----------| | `EulerDiscreteScheduler` | 20-50 | Good | Default choice | | `EulerAncestralDiscreteScheduler` | 20-50 | Good | More variation | | `DPMSolverMultistepScheduler` | 15-25 | Excellent | Fast, high quality | | `DDIMScheduler` | 50-100 | Good | Deterministic | | `LCMScheduler` | 4-8 | Good | Very fast | | `UniPCMultistepScheduler` | 15-25 | Excellent | Fast convergence | ### Swapping schedulers ```python from diffusers import DPMSolverMultistepScheduler # Swap for faster generation pipe.scheduler = DPMSolverMultistepScheduler.from_config( pipe.scheduler.config ) # Now generate with fewer steps image = pipe(prompt, num_inference_steps=20).images[0] ``` ## Generation parameters ### Key parameters | Parameter | Default | Description | |-----------|---------|-------------| | `prompt` | Required | Text description of desired image | | `negative_prompt` | None | What to avoid in the image | | `num_inference_steps` | 50 | Denoising steps (more = better quality) | | `guidance_scale` | 7.5 | Prompt adherence (7-12 typical) | | `height`, `width` | 512/1024 | Output dimensions (multiples of 8) | | `generator` | None | Torch generator for reproducibility | | `num_images_per_prompt` | 1 | Batch size | ### Reproducible generation ```python import torch generator = torch.Generator(device="cuda").manual_seed(42) image = pipe( prompt="A cat wearing a top hat", generator=generator, num_inference_steps=50 ).images[0] ``` ### Negative prompts ```python image = pipe( prompt="Professional photo of a dog in a garden", negative_prompt="blurry, low quality, distorted, ugly, bad anatomy", guidance_scale=7.5 ).images[0] ``` ## Image-to-image Transform existing images with text guidance: ```python from diffusers import AutoPipelineForImage2Image from PIL import Image pipe = AutoPipelineForImage2Image.from_pretrained( "stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16 ).to("cuda") init_image = Image.open("input.jpg").resize((512, 512)) image = pipe( prompt="A watercolor painting of the scene", image=init_image, strength=0.75, # How much to transform (0-1) num_inference_steps=50 ).images[0] ``` ## Inpainting Fill masked regions: ```python from diffusers import AutoPipelineForInpainting from PIL import Image pipe = AutoPipelineForInpainting.from_pretrained( "runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16 ).to("cuda") image = Image.open("photo.jpg") mask = Image.open("mask.png") # White = inpaint region result = pipe( prompt="A red car parked on the street", image=image, mask_image=mask, num_inference_steps=50 ).images[0] ``` ## ControlNet Add spatial conditioning for precise control: ```python from diffusers import StableDiffusionControlNetPipeline, ControlNetModel import torch # Load ControlNet for edge conditioning controlnet = ControlNetModel.from_pretrained( "lllyasviel/control_v11p_sd15_canny", torch_dtype=torch.float16 ) pipe = StableDiffusionControlNetPipeline.from_pretrained( "stable-diffusion-v1-5/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16 ).to("cuda") # Use Canny edge image as control control_image = get_canny_image(input_image) image = pipe( prompt="A beautiful house in the style of Van Gogh", image=control_image, num_inference_steps=30 ).images[0] ``` ### Available ControlNets | ControlNet | Input Type | Use Case | |------------|------------|----------| | `canny` | Edge maps | Preserve structure | | `openpose` | Pose skeletons | Human poses | | `depth` | Depth maps | 3D-aware generation | | `normal` | Normal maps | Surface details | | `mlsd` | Line segments | Architectural lines | | `scribble` | Rough sketches | Sketch-to-image | ## LoRA adapters Load fine-tuned style adapters: ```python from diffusers import DiffusionPipeline pipe = DiffusionPipeline.from_pretrained( "stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16 ).to("cuda") # Load LoRA weights pipe.load_lora_weights("path/to/lora", weight_name="style.safetensors") # Generate with LoRA style image = pipe("A portrait in the trained style").images[0] # Adjust LoRA strength pipe.fuse_lora(lora_scale=0.8) # Unload LoRA pipe.unload_lora_weights() ``` ### Multiple LoRAs ```python # Load multiple LoRAs pipe.load_lora_weights("lora1", adapter_name="style") pipe.load_lora_weights("lora2", adapter_name="character") # Set weights for each pipe.set_adapters(["style", "character"], adapter_weights=[0.7, 0.5]) image = pipe("A portrait").images[0] ``` ## Memory optimization ### Enable CPU offloading ```python # Model CPU offload - moves models to CPU when not in use pipe.enable_model_cpu_offload() # Sequential CPU offload - more aggressive, slower pipe.enable_sequential_cpu_offload() ``` ### Attention slicing ```python # Reduce memory by computing attention in chunks pipe.enable_attention_slicing() # Or specific chunk size pipe.enable_attention_slicing("max") ``` ### xFormers memory-efficient attention ```python # Requires xformers package pipe.enable_xformers_memory_efficient_attention() ``` ### VAE slicing for large images ```python # Decode latents in tiles for large images pipe.enable_vae_slicing() pipe.enable_vae_tiling() ``` ## Model variants ### Loading different precisions ```python # FP16 (recommended for GPU) pipe = DiffusionPipeline.from_pretrained( "model-id", torch_dtype=torch.float16, variant="fp16" ) # BF16 (better precision, requires Ampere+ GPU) pipe = DiffusionPipeline.from_pretrained( "model-id", torch_dtype=torch.bfloat16 ) ``` ### Loading specific components ```python from diffusers import UNet2DConditionModel, AutoencoderKL # Load custom VAE vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse") # Use with pipeline pipe = DiffusionPipeline.from_pretrained( "stable-diffusion-v1-5/stable-diffusion-v1-5", vae=vae, torch_dtype=torch.float16 ) ``` ## Batch generation Generate multiple images efficiently: ```python # Multiple prompts prompts = [ "A cat playing piano", "A dog reading a book", "A bird painting a picture" ] images = pipe(prompts, num_inference_steps=30).images # Multiple images per prompt images = pipe( "A beautiful sunset", num_images_per_prompt=4, num_inference_steps=30 ).images ``` ## Common workflows ### Workflow 1: High-quality generation ```python from diffusers import StableDiffusionXLPipeline, DPMSolverMultistepScheduler import torch # 1. Load SDXL with optimizations pipe = StableDiffusionXLPipeline.from_pretrained( "stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16" ) pipe.to("cuda") pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config) pipe.enable_model_cpu_offload() # 2. Generate with quality settings image = pipe( prompt="A majestic lion in the savanna, golden hour lighting, 8k, detailed fur", negative_prompt="blurry, low quality, cartoon, anime, sketch", num_inference_steps=30, guidance_scale=7.5, height=1024, width=1024 ).images[0] ``` ### Workflow 2: Fast prototyping ```python from diffusers import AutoPipelineForText2Image, LCMScheduler import torch # Use LCM for 4-8 step generation pipe = AutoPipelineForText2Image.from_pretrained( "stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16 ).to("cuda") # Load LCM LoRA for fast generation pipe.load_lora_weights("latent-consistency/lcm-lora-sdxl") pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config) pipe.fuse_lora() # Generate in ~1 second image = pipe( "A beautiful landscape", num_inference_steps=4, guidance_scale=1.0 ).images[0] ``` ## Common issues **CUDA out of memory:** ```python # Enable memory optimizations pipe.enable_model_cpu_offload() pipe.enable_attention_slicing() pipe.enable_vae_slicing() # Or use lower precision